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from typing import Any, Callable, Dict, List, Optional, Tuple, Union |
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import torch |
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from diffusers import StableDiffusionXLPipeline |
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from diffusers.pipelines.stable_diffusion_xl import StableDiffusionXLPipelineOutput |
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from diffusers.pipelines.stable_diffusion_xl.pipeline_stable_diffusion_xl import rescale_noise_cfg |
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from .utils import is_torch2_available |
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if is_torch2_available(): |
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from .attention_processor import IPAttnProcessor2_0 as IPAttnProcessor |
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else: |
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from .attention_processor import IPAttnProcessor |
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class StableDiffusionXLCustomPipeline(StableDiffusionXLPipeline): |
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def set_scale(self, scale): |
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for attn_processor in self.unet.attn_processors.values(): |
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if isinstance(attn_processor, IPAttnProcessor): |
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attn_processor.scale = scale |
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@torch.no_grad() |
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def __call__( |
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self, |
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prompt: Optional[Union[str, List[str]]] = None, |
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prompt_2: Optional[Union[str, List[str]]] = None, |
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height: Optional[int] = None, |
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width: Optional[int] = None, |
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num_inference_steps: int = 50, |
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denoising_end: Optional[float] = None, |
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guidance_scale: float = 5.0, |
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negative_prompt: Optional[Union[str, List[str]]] = None, |
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negative_prompt_2: Optional[Union[str, List[str]]] = None, |
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num_images_per_prompt: Optional[int] = 1, |
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eta: float = 0.0, |
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generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None, |
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latents: Optional[torch.FloatTensor] = None, |
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prompt_embeds: Optional[torch.FloatTensor] = None, |
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negative_prompt_embeds: Optional[torch.FloatTensor] = None, |
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pooled_prompt_embeds: Optional[torch.FloatTensor] = None, |
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negative_pooled_prompt_embeds: Optional[torch.FloatTensor] = None, |
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output_type: Optional[str] = "pil", |
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return_dict: bool = True, |
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callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None, |
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callback_steps: int = 1, |
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cross_attention_kwargs: Optional[Dict[str, Any]] = None, |
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guidance_rescale: float = 0.0, |
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original_size: Optional[Tuple[int, int]] = None, |
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crops_coords_top_left: Tuple[int, int] = (0, 0), |
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target_size: Optional[Tuple[int, int]] = None, |
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negative_original_size: Optional[Tuple[int, int]] = None, |
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negative_crops_coords_top_left: Tuple[int, int] = (0, 0), |
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negative_target_size: Optional[Tuple[int, int]] = None, |
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control_guidance_start: float = 0.0, |
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control_guidance_end: float = 1.0, |
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): |
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r""" |
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Function invoked when calling the pipeline for generation. |
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Args: |
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prompt (`str` or `List[str]`, *optional*): |
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The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`. |
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instead. |
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prompt_2 (`str` or `List[str]`, *optional*): |
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The prompt or prompts to be sent to the `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is |
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used in both text-encoders |
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height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor): |
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The height in pixels of the generated image. This is set to 1024 by default for the best results. |
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Anything below 512 pixels won't work well for |
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[stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) |
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and checkpoints that are not specifically fine-tuned on low resolutions. |
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width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor): |
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The width in pixels of the generated image. This is set to 1024 by default for the best results. |
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Anything below 512 pixels won't work well for |
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[stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) |
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and checkpoints that are not specifically fine-tuned on low resolutions. |
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num_inference_steps (`int`, *optional*, defaults to 50): |
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The number of denoising steps. More denoising steps usually lead to a higher quality image at the |
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expense of slower inference. |
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denoising_end (`float`, *optional*): |
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When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be |
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completed before it is intentionally prematurely terminated. As a result, the returned sample will |
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still retain a substantial amount of noise as determined by the discrete timesteps selected by the |
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scheduler. The denoising_end parameter should ideally be utilized when this pipeline forms a part of a |
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"Mixture of Denoisers" multi-pipeline setup, as elaborated in [**Refining the Image |
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Output**](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/stable_diffusion_xl#refining-the-image-output) |
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guidance_scale (`float`, *optional*, defaults to 5.0): |
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Guidance scale as defined in [Classifier-Free Diffusion Guidance](https://arxiv.org/abs/2207.12598). |
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`guidance_scale` is defined as `w` of equation 2. of [Imagen |
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Paper](https://arxiv.org/pdf/2205.11487.pdf). Guidance scale is enabled by setting `guidance_scale > |
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1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`, |
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usually at the expense of lower image quality. |
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negative_prompt (`str` or `List[str]`, *optional*): |
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The prompt or prompts not to guide the image generation. If not defined, one has to pass |
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`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is |
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less than `1`). |
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negative_prompt_2 (`str` or `List[str]`, *optional*): |
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The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and |
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`text_encoder_2`. If not defined, `negative_prompt` is used in both text-encoders |
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num_images_per_prompt (`int`, *optional*, defaults to 1): |
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The number of images to generate per prompt. |
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eta (`float`, *optional*, defaults to 0.0): |
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Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to |
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[`schedulers.DDIMScheduler`], will be ignored for others. |
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generator (`torch.Generator` or `List[torch.Generator]`, *optional*): |
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One or a list of [torch generator(s)](https://pytorch.org/docs/stable/generated/torch.Generator.html) |
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to make generation deterministic. |
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latents (`torch.FloatTensor`, *optional*): |
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Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image |
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generation. Can be used to tweak the same generation with different prompts. If not provided, a latents |
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tensor will ge generated by sampling using the supplied random `generator`. |
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prompt_embeds (`torch.FloatTensor`, *optional*): |
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Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not |
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provided, text embeddings will be generated from `prompt` input argument. |
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negative_prompt_embeds (`torch.FloatTensor`, *optional*): |
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Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt |
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weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input |
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argument. |
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pooled_prompt_embeds (`torch.FloatTensor`, *optional*): |
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Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. |
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If not provided, pooled text embeddings will be generated from `prompt` input argument. |
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negative_pooled_prompt_embeds (`torch.FloatTensor`, *optional*): |
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Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt |
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weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` |
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input argument. |
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output_type (`str`, *optional*, defaults to `"pil"`): |
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The output format of the generate image. Choose between |
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[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`. |
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return_dict (`bool`, *optional*, defaults to `True`): |
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Whether or not to return a [`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput`] instead |
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of a plain tuple. |
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callback (`Callable`, *optional*): |
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A function that will be called every `callback_steps` steps during inference. The function will be |
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called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`. |
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callback_steps (`int`, *optional*, defaults to 1): |
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The frequency at which the `callback` function will be called. If not specified, the callback will be |
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called at every step. |
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cross_attention_kwargs (`dict`, *optional*): |
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A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under |
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`self.processor` in |
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[diffusers.models.attention_processor](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py). |
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guidance_rescale (`float`, *optional*, defaults to 0.7): |
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Guidance rescale factor proposed by [Common Diffusion Noise Schedules and Sample Steps are |
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Flawed](https://arxiv.org/pdf/2305.08891.pdf) `guidance_scale` is defined as `φ` in equation 16. of |
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[Common Diffusion Noise Schedules and Sample Steps are Flawed](https://arxiv.org/pdf/2305.08891.pdf). |
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Guidance rescale factor should fix overexposure when using zero terminal SNR. |
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original_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): |
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If `original_size` is not the same as `target_size` the image will appear to be down- or upsampled. |
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`original_size` defaults to `(width, height)` if not specified. Part of SDXL's micro-conditioning as |
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explained in section 2.2 of |
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[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). |
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crops_coords_top_left (`Tuple[int]`, *optional*, defaults to (0, 0)): |
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`crops_coords_top_left` can be used to generate an image that appears to be "cropped" from the position |
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`crops_coords_top_left` downwards. Favorable, well-centered images are usually achieved by setting |
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`crops_coords_top_left` to (0, 0). Part of SDXL's micro-conditioning as explained in section 2.2 of |
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[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). |
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target_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): |
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For most cases, `target_size` should be set to the desired height and width of the generated image. If |
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not specified it will default to `(width, height)`. Part of SDXL's micro-conditioning as explained in |
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section 2.2 of [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). |
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negative_original_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): |
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To negatively condition the generation process based on a specific image resolution. Part of SDXL's |
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micro-conditioning as explained in section 2.2 of |
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[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more |
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information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. |
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negative_crops_coords_top_left (`Tuple[int]`, *optional*, defaults to (0, 0)): |
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To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's |
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micro-conditioning as explained in section 2.2 of |
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[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more |
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information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. |
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negative_target_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): |
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To negatively condition the generation process based on a target image resolution. It should be as same |
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as the `target_size` for most cases. Part of SDXL's micro-conditioning as explained in section 2.2 of |
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[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more |
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information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. |
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control_guidance_start (`float`, *optional*, defaults to 0.0): |
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The percentage of total steps at which the ControlNet starts applying. |
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control_guidance_end (`float`, *optional*, defaults to 1.0): |
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The percentage of total steps at which the ControlNet stops applying. |
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Examples: |
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Returns: |
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[`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput`] or `tuple`: |
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[`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput`] if `return_dict` is True, otherwise a |
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`tuple`. When returning a tuple, the first element is a list with the generated images. |
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""" |
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height = height or self.default_sample_size * self.vae_scale_factor |
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width = width or self.default_sample_size * self.vae_scale_factor |
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original_size = original_size or (height, width) |
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target_size = target_size or (height, width) |
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self.check_inputs( |
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prompt, |
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prompt_2, |
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height, |
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width, |
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callback_steps, |
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negative_prompt, |
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negative_prompt_2, |
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prompt_embeds, |
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negative_prompt_embeds, |
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pooled_prompt_embeds, |
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negative_pooled_prompt_embeds, |
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) |
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if prompt is not None and isinstance(prompt, str): |
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batch_size = 1 |
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elif prompt is not None and isinstance(prompt, list): |
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batch_size = len(prompt) |
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else: |
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batch_size = prompt_embeds.shape[0] |
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device = self._execution_device |
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do_classifier_free_guidance = guidance_scale > 1.0 |
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text_encoder_lora_scale = ( |
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cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None |
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) |
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( |
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prompt_embeds, |
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negative_prompt_embeds, |
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pooled_prompt_embeds, |
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negative_pooled_prompt_embeds, |
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) = self.encode_prompt( |
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prompt=prompt, |
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prompt_2=prompt_2, |
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device=device, |
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num_images_per_prompt=num_images_per_prompt, |
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do_classifier_free_guidance=do_classifier_free_guidance, |
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negative_prompt=negative_prompt, |
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negative_prompt_2=negative_prompt_2, |
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prompt_embeds=prompt_embeds, |
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negative_prompt_embeds=negative_prompt_embeds, |
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pooled_prompt_embeds=pooled_prompt_embeds, |
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negative_pooled_prompt_embeds=negative_pooled_prompt_embeds, |
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lora_scale=text_encoder_lora_scale, |
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) |
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self.scheduler.set_timesteps(num_inference_steps, device=device) |
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timesteps = self.scheduler.timesteps |
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num_channels_latents = self.unet.config.in_channels |
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latents = self.prepare_latents( |
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batch_size * num_images_per_prompt, |
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num_channels_latents, |
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height, |
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width, |
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prompt_embeds.dtype, |
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device, |
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generator, |
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latents, |
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) |
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extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta) |
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add_text_embeds = pooled_prompt_embeds |
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if self.text_encoder_2 is None: |
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text_encoder_projection_dim = int(pooled_prompt_embeds.shape[-1]) |
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else: |
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text_encoder_projection_dim = self.text_encoder_2.config.projection_dim |
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add_time_ids = self._get_add_time_ids( |
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original_size, |
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crops_coords_top_left, |
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target_size, |
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dtype=prompt_embeds.dtype, |
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text_encoder_projection_dim=text_encoder_projection_dim, |
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) |
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if negative_original_size is not None and negative_target_size is not None: |
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negative_add_time_ids = self._get_add_time_ids( |
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negative_original_size, |
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negative_crops_coords_top_left, |
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negative_target_size, |
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dtype=prompt_embeds.dtype, |
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text_encoder_projection_dim=text_encoder_projection_dim, |
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) |
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else: |
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negative_add_time_ids = add_time_ids |
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if do_classifier_free_guidance: |
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prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds], dim=0) |
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add_text_embeds = torch.cat([negative_pooled_prompt_embeds, add_text_embeds], dim=0) |
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add_time_ids = torch.cat([negative_add_time_ids, add_time_ids], dim=0) |
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prompt_embeds = prompt_embeds.to(device) |
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add_text_embeds = add_text_embeds.to(device) |
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add_time_ids = add_time_ids.to(device).repeat(batch_size * num_images_per_prompt, 1) |
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num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0) |
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if denoising_end is not None and isinstance(denoising_end, float) and denoising_end > 0 and denoising_end < 1: |
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discrete_timestep_cutoff = int( |
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round( |
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self.scheduler.config.num_train_timesteps |
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- (denoising_end * self.scheduler.config.num_train_timesteps) |
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) |
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) |
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num_inference_steps = len(list(filter(lambda ts: ts >= discrete_timestep_cutoff, timesteps))) |
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timesteps = timesteps[:num_inference_steps] |
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for attn_processor in self.unet.attn_processors.values(): |
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if isinstance(attn_processor, IPAttnProcessor): |
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conditioning_scale = attn_processor.scale |
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break |
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with self.progress_bar(total=num_inference_steps) as progress_bar: |
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for i, t in enumerate(timesteps): |
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if (i / len(timesteps) < control_guidance_start) or ((i + 1) / len(timesteps) > control_guidance_end): |
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self.set_scale(0.0) |
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else: |
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self.set_scale(conditioning_scale) |
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latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents |
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latent_model_input = self.scheduler.scale_model_input(latent_model_input, t) |
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added_cond_kwargs = {"text_embeds": add_text_embeds, "time_ids": add_time_ids} |
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noise_pred = self.unet( |
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latent_model_input, |
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t, |
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encoder_hidden_states=prompt_embeds, |
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cross_attention_kwargs=cross_attention_kwargs, |
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added_cond_kwargs=added_cond_kwargs, |
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return_dict=False, |
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)[0] |
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if do_classifier_free_guidance: |
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noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) |
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noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) |
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if do_classifier_free_guidance and guidance_rescale > 0.0: |
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noise_pred = rescale_noise_cfg(noise_pred, noise_pred_text, guidance_rescale=guidance_rescale) |
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latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0] |
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if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0): |
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progress_bar.update() |
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if callback is not None and i % callback_steps == 0: |
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callback(i, t, latents) |
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|
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if not output_type == "latent": |
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|
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needs_upcasting = self.vae.dtype == torch.float16 and self.vae.config.force_upcast |
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|
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if needs_upcasting: |
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self.upcast_vae() |
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latents = latents.to(next(iter(self.vae.post_quant_conv.parameters())).dtype) |
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image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0] |
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if needs_upcasting: |
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self.vae.to(dtype=torch.float16) |
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else: |
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image = latents |
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if output_type != "latent": |
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if self.watermark is not None: |
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image = self.watermark.apply_watermark(image) |
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image = self.image_processor.postprocess(image, output_type=output_type) |
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self.maybe_free_model_hooks() |
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if not return_dict: |
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return (image,) |
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return StableDiffusionXLPipelineOutput(images=image) |
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