SPO-Diffusion-Models/SPO-SDXL_4k-p_10ep_LoRA

Step-aware Preference Optimization: Aligning Preference with Denoising Performance at Each Step

Abstract

Recently, Direct Preference Optimization (DPO) has extended its success from aligning large language models (LLMs) to aligning text-to-image diffusion models with human preferences. Unlike most existing DPO methods that assume all diffusion steps share a consistent preference order with the final generated images, we argue that this assumption neglects step-specific denoising performance and that preference labels should be tailored to each step's contribution.

To address this limitation, we propose Step-aware Preference Optimization (SPO), a novel post-training approach that independently evaluates and adjusts the denoising performance at each step, using a step-aware preference model and a step-wise resampler to ensure accurate step-aware supervision. Specifically, at each denoising step, we sample a pool of images, find a suitable win-lose pair, and, most importantly, randomly select a single image from the pool to initialize the next denoising step. This step-wise resampler process ensures the next win-lose image pair comes from the same image, making the win-lose comparison independent of the previous step. To assess the preferences at each step, we train a separate step-aware preference model that can be applied to both noisy and clean images.

Our experiments with Stable Diffusion v1.5 and SDXL demonstrate that SPO significantly outperforms the latest Diffusion-DPO in aligning generated images with complex, detailed prompts and enhancing aesthetics, while also achieving more than 20× times faster in training efficiency. Code and model: https://rockeycoss.github.io/spo.github.io/

Model Description

This model is fine-tuned from stable-diffusion-xl-base-1.0. It has been trained on 4,000 prompts for 10 epochs. This checkpoint is a LoRA checkpoint. We also provide a LoRA checkpoint compatible with stable-diffusion-webui, which can be accessed here

If you want to access the merged checkpoint that combines the LoRA checkpoint with the base model stable-diffusion-xl-base-1.0, please visit SPO-SDXL_4k-p_10ep.

A quick example

from diffusers import StableDiffusionXLPipeline, UNet2DConditionModel
import torch

# load pipeline
inference_dtype = torch.float16
pipe = StableDiffusionXLPipeline.from_pretrained(
    "stabilityai/stable-diffusion-xl-base-1.0",
    torch_dtype=inference_dtype,
)
vae = AutoencoderKL.from_pretrained(
    'madebyollin/sdxl-vae-fp16-fix',
    torch_dtype=inference_dtype,
)
pipe.vae = vae
pipe.load_lora_weights("SPO-Diffusion-Models/SPO-SDXL_4k-p_10ep_LoRA")
pipe.to('cuda')

generator=torch.Generator(device='cuda').manual_seed(42)
image = pipe(
    prompt='a child and a penguin sitting in front of the moon',
    guidance_scale=5.0,
    generator=generator,
    output_type='pil',
).images[0]
image.save('moon.png')

Citation

If you find our work or codebase useful, please consider giving us a star and citing our work.

@article{liang2024step,
  title={Step-aware Preference Optimization: Aligning Preference with Denoising Performance at Each Step},
  author={Liang, Zhanhao and Yuan, Yuhui and Gu, Shuyang and Chen, Bohan and Hang, Tiankai and Li, Ji and Zheng, Liang},
  journal={arXiv preprint arXiv:2406.04314},
  year={2024}
}